K = r × u = 4 × 5 = 20.Now, u = 1/6, substitute this value in the above equation.r = k/u = 20/(1/6) = 120, if u = 1/6, then r = 120.
Given that r varies inversely as u and r = 4 when u = 5. To find the value of r when u = 1/6. Inversely proportional variables: When one variable increases and the other variable decreases, then two variables are said to be inversely proportional to each other. It can be shown as:r α 1/u ⇒ r = k/uwhere k is the constant of variation. Here, k = r × u. We know that when u = 5, r = 4. Therefore, k = r × u = 4 × 5 = 20.Now, u = 1/6, substitute this value in the above equation.r = k/u = 20/(1/6) = 120Hence, the value of r is 120 when u = 1/6.Answer:Therefore, if u = 1/6, then r = 120.
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Find the area of a triangle with sides 7 yards, 7 yards, and 5 yards. (Round your answer to one decimal place.)
The area of the triangle with sides 7 yards, 7 yards, and 5 yards is approximately 17.1 square yards. To find the area of a triangle, we can use Heron's formula, which states that the area (A) of a triangle with sides a, b, and c can be calculated using the semi-perimeter (s) of the triangle.
The semi-perimeter of a triangle is:
s = (a + b + c) / 2
The area can then be calculated as:
A = √(s(s - a)(s - b)(s - c))
Given the sides of the triangle as 7 yards, 7 yards, and 5 yards, we can calculate the semi-perimeter:
s = (7 + 7 + 5) / 2
s = 19 / 2
s = 9.5 yards
Using this value, we can calculate the area:
A = √(9.5(9.5 - 7)(9.5 - 7)(9.5 - 5))
A = √(9.5 * 2.5 * 2.5 * 4.5)
A ≈ √(237.1875)
A ≈ 15.4 square yards
Rounding this value to one decimal place, the area of the triangle is approximately 17.1 square yards.
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red n Let Ao be an 4 x 4-matrix with det (Ao) = 3. Compute the determinant of the matrices A1, A2, A3, A4 and A5, obtained from Ao by the following operations: A₁ is obtained from Ao by multiplying the fourth row of Ao by the number 3. det (A₁) = [2mark] A2 is obtained from Ao by replacing the second row by the sum of itself plus the 4 times the third row. det (A₂) = [2mark] A3 is obtained from Ao by multiplying Ao by itself.. det (A3) = [2mark] A4 is obtained from Ao by swapping the first and last rows of Ag. A2 is obtained from Ao by replacing the second row by the sum of itself plus the 4 times the third row. det (A₂) = [2mark] A3 is obtained from Ao by multiplying Ao by itself.. det (A3) = [2mark] A4 is obtained from Ao by swapping the first and last rows of Ag. det (A4) = [2mark] As is obtained from Ao by scaling Ao by the number 2. det (A5) = [2mark]
Given a 4x4 matrix [tex]A_{o}[/tex] with det([tex]A_{o}[/tex]) = 3, we need to compute the determinants of the matrices [tex]A_{1}[/tex], [tex]A_{2}[/tex], [tex]A_{3[/tex], [tex]A_{4}[/tex], and [tex]A_{5}[/tex], obtained by performing specific operations on [tex]A_{o}[/tex].
The determinants are as follows: det([tex]A_{1}[/tex]) = ?, det([tex]A_{2}[/tex]) = ?, det([tex]A_{3[/tex]) = ?, det( [tex]A_{4}[/tex]) = ?, det([tex]A_{5}[/tex]}) = ?
To compute the determinants of the matrices obtained from [tex]A_{o}[/tex] by different operations, let's go through each operation:
[tex]A_{1}[/tex] is obtained by multiplying the fourth row of [tex]A_{o}[/tex] by 3:
To find det([tex]A_{1}[/tex]), we can simply multiply the determinant of [tex]A_{o}[/tex] by 3 since multiplying a row by a scalar multiplies the determinant by the same scalar. Therefore, det([tex]A_{1}[/tex]) = 3 * det([tex]A_{o}[/tex]) = 3 * 3 = 9.
[tex]A_{2}[/tex] is obtained by replacing the second row with the sum of itself and 4 times the third row:
This operation does not affect the determinant since adding a multiple of one row to another does not change the determinant. Hence, det([tex]A_{2}[/tex]) = det([tex]A_{o}[/tex]) = 3.
[tex]A_{3[/tex] is obtained by multiplying [tex]A_{o}[/tex] by itself:
When multiplying two matrices, the determinant of the resulting matrix is the product of the determinants of the original matrices. Thus, det([tex]A_{3[/tex]) = det([tex]A_{o}[/tex]) * det([tex]A_{o}[/tex]) = 3 * 3 = 9.
[tex]A_{4}[/tex] is obtained by swapping the first and last rows of [tex]A_{o}[/tex]:
Swapping rows changes the sign of the determinant, so det([tex]A_{4}[/tex]) = -det([tex]A_{o}[/tex]) = -3.
[tex]A_{5}[/tex] is obtained by scaling [tex]A_{o}[/tex] by 2:
Similar to [tex]A_{1}[/tex], scaling a row multiplies the determinant by the same scalar. Therefore, det([tex]A_{5}[/tex]) = 2 * det([tex]A_{o}[/tex]) = 2 * 3 = 6.
In summary, the determinants of the matrices are: det([tex]A_{1}[/tex]) = 9, det([tex]A_{2}[/tex]) = 3, det([tex]A_{3[/tex]) = 9, det( [tex]A_{4}[/tex]) = -3, and det([tex]A_{5}[/tex]) = 6.
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find the absolute maximum and minimum values of the function over the indicated interval, and indicate the x-values at which they occur f(x)=x^2-4x-9; [0,5]
The absolute maximum and minimum values of the function over the indicated interval and indicate the x-values at which they occur f(x) = x² - 4x - 9; [0, 5],
we need to follow the steps given below:
Step 1: Differentiate the given function to find the critical points and intervals where the function increases and decreases.
f(x) = x² - 4x - 9f'(x)
= 2x - 4= 0
⇒ 2x = 4
⇒ x = 2
Thus, we get a critical point at x = 2.
Now, we will find the intervals where the function increases and decreases using the test point method:
f'(x) = 2x - 4> 0 for x > 2
∴ f(x) is increasing for x > 2.f'(x) = 2x - 4< 0 for x < 2
∴ f(x) is decreasing for x < 2.
Step 2: Check the function values at the critical points and the end points of the interval.
f(0) = (0)² - 4(0) - 9
= -9f(2) = (2)² - 4(2) - 9
= -13f(5) = (5)² - 4(5) - 9
= -19
Step 3: Now, we can identify the absolute maximum and minimum values of the function over the indicated interval
[0, 5].
Absolute maximum value of the function:
The absolute maximum value of the function over the interval [0, 5] is -9 and it occurs at x = 0.
Absolute minimum value of the function:
The absolute minimum value of the function over the interval [0, 5] is -19 and it occurs at x = 5.
Therefore, the absolute maximum and minimum values of the function over the indicated interval [0, 5] and the x-values at which they occur are as follows.
Absolute maximum value = -9 at x = 0
Absolute minimum value = -19 at x = 5
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Find the circumference of the circle.
radius is 12cm
Circumference of circle is,
⇒ C = 75.36 cm
We have to given that,
Radius of circle is,
⇒ r = 12 cm
Since, We know that,
Circumference of circle is,
⇒ C = 2πr
Where, 'r' is radius and π is 3.14,
Here, we have;
⇒ r = 12 cm
Hence, We get;
Circumference of circle is,
⇒ C = 2πr
⇒ C = 2 × 3.14 × 12
⇒ C = 75.36 cm
Therefore, Circumference of circle is,
⇒ C = 2πr
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Not yet answered Marked out of 1.00 Question 8 Let A and B be events in a random experiment. Suppose that A and B are independent and P(A) = 0.4 and P(B) = 0.2. Then P(A - B) = Select one: none a. b. 0.32 0.18 C. d. 0.12
A and B be events in a random experiment. The correct answer is (b) 0.32.
To find P(A - B), we need to subtract the probability of event B from the probability of event A. In other words, we want to find the probability of event A occurring without the occurrence of event B.
Since A and B are independent events, the probability of their intersection (A ∩ B) is equal to the product of their individual probabilities: P(A ∩ B) = P(A) * P(B).
We can use this information to find P(A - B) as follows:
P(A - B) = P(A) - P(A ∩ B)
Since A and B are independent, P(A ∩ B) = P(A) * P(B).
P(A - B) = P(A) - P(A) * P(B)
Given that P(A) = 0.4 and P(B) = 0.2, we can substitute these values into the equation:
P(A - B) = 0.4 - 0.4 * 0.2
P(A - B) = 0.4 - 0.08
P(A - B) = 0.32
Therefore, the correct answer is (b) 0.32.
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are and homogeneous coordinates for the same point in ? why or why not?
No, Euclidean coordinates and homogeneous coordinates are not the same thing for the same point in space. Let's see how are they different in this brief discussion below. What are homogeneous coordinates? Homogeneous coordinates are utilized to explain geometry in projective space. Homogeneous coordinates are often used since they can express points at infinity. Homogeneous coordinates are three-dimensional coordinates used to extend projective space to include points at infinity. How are homogeneous coordinates and Euclidean coordinates different?Homogeneous coordinates utilize four variables to define a point in space while Euclidean coordinates use three variables. Points in Euclidean geometry have no "weights" or "scales," while points in projective geometry can be "scaled" to make them homogeneous. Hence, Euclidean coordinates and homogeneous coordinates are not the same thing for the same point in space.
Homogeneous coordinates and Cartesian coordinates are not the same point.
The following are the reasons behind it:
Homogeneous coordinates :Homogeneous coordinates are a set of coordinates in which the value of any point in space is represented by three coordinates in a ratio, which means that the first two coordinates can be increased or decreased in size, but the third coordinate should also be changed proportionally.
So, in short, these are different representations of the same point. Homogeneous coordinates are used in 3D modeling, computer vision, and other applications.
Cartesian coordinates: Cartesian coordinates, also known as rectangular coordinates, are the usual (x, y) coordinates.
These coordinates are widely used in mathematics to explain the relationship between geometric shapes and points. These are the coordinate points that we use in our daily lives, such as identifying the location of a particular spot on a map or finding the shortest path between two points on a coordinate plane.
The two-dimensional (2D) or three-dimensional (3D) points are represented by Cartesian coordinates.
Hence, it can be concluded that Homogeneous coordinates and Cartesian coordinates are not the same point, and these are different representations of the same point.
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Suppose that we have 100 apples. In order to determine the integrity of the entire batch of apples, we carefully examine n randomly-chosen apples; if any of the apples is rotten, the whole batch of apples is discarded. Suppose that 50 of the apples are rotten, but we do not know this during the inspection process.
(a) Calculate the probability that the whole batch is discarded for n = 1, 2, 3, 4, 5, 6
(b) Find all values of n for which the probability of discarding the whole batch of apples is at least 99% = 99/100
(a) To calculate the probability that the whole batch is discarded for a given value of n, we need to consider the probability that at least one of the randomly chosen apples is rotten.
Let's calculate this probability for each value of n:
For n = 1:
The probability that at least one apple is rotten is 50/100 = 1/2.
Therefore, the probability that the whole batch is discarded is 1/2.
For n = 2:
The probability that both apples are not rotten is (50/100) * (49/99) = 2450/9900.
Therefore, the probability that at least one apple is rotten is 1 - (2450/9900) = 7450/9900.
Therefore, the probability that the whole batch is discarded is 7450/9900.
For n = 3:
The probability that all three apples are not rotten is (50/100) * (49/99) * (48/98) = 117600/485100.
Therefore, the probability that at least one apple is rotten is 1 - (117600/485100) = 367500/485100.
Therefore, the probability that the whole batch is discarded is 367500/485100.
For n = 4:
The probability that all four apples are not rotten is (50/100) * (49/99) * (48/98) * (47/97) = 342200/1088433.
Therefore, the probability that at least one apple is rotten is 1 - (342200/1088433) = 746233/1088433.
Therefore, the probability that the whole batch is discarded is 746233/1088433.
For n = 5:
The probability that all five apples are not rotten is (50/100) * (49/99) * (48/98) * (47/97) * (46/96) = 50702400/182530530.
Therefore, the probability that at least one apple is rotten is 1 - (50702400/182530530) = 131828130/182530530.
Therefore, the probability that the whole batch is discarded is 131828130/182530530.
For n = 6:
The probability that all six apples are not rotten is (50/100) * (49/99) * (48/98) * (47/97) * (46/96) * (45/95) = 386914800/1251677705.
Therefore, the probability that at least one apple is rotten is 1 - (386914800/1251677705) = 864762905/1251677705.
Therefore, the probability that the whole batch is discarded is 864762905/1251677705.
(b) To find the values of n for which the probability of discarding the whole batch of apples is at least 99/100, we need to find the smallest value of n such that the probability exceeds or equals 99/100.
Starting from n = 1, we can calculate the probability for each value of n until we reach a probability greater than or equal to 99/100:
For n = 1: Probability = 1/2.
For n = 2: Probability = 7450/9900.
For n = 3: Probability = 367500/485100.
For n = 4: Probability = 746233/1088433.
For n = 5: Probability = 131828130/182530530.
For n = 6: Probability = 864762905/1251677705.
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A medical researcher believes that the variance of total cholesterol levels in men is greater than the variance of total cholesterol levels in women. The sample variance for a random sample of 9 men’s cholesterol levels, measured in mgdL, is 287. The sample variance for a random sample of 8 women is 88. Assume that both population distributions are approximately normal and test the researcher’s claim using a 0.10 level of significance. Does the evidence support the researcher’s belief? Let men's total cholesterol levels be Population 1 and let women's total cholesterol levels be Population 2.
1 State the null and alternative hypotheses for the test. Fill in the blank below. H0Ha: σ21=σ22: σ21⎯⎯⎯⎯⎯⎯⎯⎯⎯⎯⎯⎯⎯⎯⎯σ22
2. What is the test statistic?
3. Draw a conclusion
The null and alternative hypotheses for the test are as follows: Null hypothesis (H 0): The variance of total cholesterol levels in men is equal to the variance of total cholesterol levels in women.
Alternative hypothesis (H a): The variance of total cholesterol levels in men is greater than the variance of total cholesterol levels in women.
The null hypothesis states that the variances of total cholesterol levels in men and women are equal, while the alternative hypothesis suggests that the variance in men is greater than that in women. The notation σ21 represents the variance of men's total cholesterol levels, and σ22 represents the variance of women's total cholesterol levels.
The test statistic for comparing variances is the F statistic, calculated as the ratio of the sample variances: F = (sample variance of men) / (sample variance of women). In this case, the sample variance of men is 287 and the sample variance of women is 88.
To draw a conclusion, we compare the calculated F statistic with the critical value from the F distribution at a significance level of 0.10. If the calculated F statistic is greater than the critical value, we reject the null hypothesis and conclude that there is evidence to support the researcher's belief that the variance of total cholesterol levels in men is greater than in women. If the calculated F statistic is not greater than the critical value, we fail to reject the null hypothesis and do not have sufficient evidence to support the researcher's belief.
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3. Find general solution. y(4) — y" = 5e² + 3 Write clean, and clear. Show steps of calculations. Hint: use the method of undetermined coefficients for the particular solution yp.
the particular solution is yp = (-5/4)e^2 + B.To find the general solution of the differential equation y(4) - y" = 5e² + 3, we'll solve for the complementary solution and the particular solution separately.
First, let's find the complementary solution by assuming y = e^(rx) and substituting it into the equation. This yields the characteristic equation r^4 - r^2 = 0. Factoring out r^2, we get r^2(r^2 - 1) = 0. So the roots are r = 0, ±1.
The complementary solution is y_c = C₁ + C₂e^x + C₃e^(-x) + C₄e^(0), which simplifies to y_c = C₁ + C₂e^x + C₃e^(-x) + C₄.
Next, we'll find the particular solution using the method of undetermined coefficients. Since the right-hand side is a combination of exponential and constant terms, we assume a particular solution of the form yp = Ae^2 + B.
Substituting this into the differential equation, we get -4Ae^2 = 5e^2 + 3. Equating the coefficients, we have -4A = 5, which gives A = -5/4.
Thus, thethe particular solution is yp = (-5/4)e^2 + B.
Combining the complementary and particular solutions, the general solution of the differential equation is y = C₁ + C₂e^x + C₃e^(-x) + C₄ + (-5/4)e^2 + B.
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fill in the blank. A particular city had a population of 27,000 in 1930 and a population of 32,000 in 1950. Assuming that its population continues to grow exponentially at a constant rate, what population will it have in 2000? The population of the city in 2000 will be people. (Round the final answer to the nearest whole number as needed. Round all intermediate values to six decimal places as needed.)
The population of the city in 2000 will be approximately 38,534 people.
How many people will be living in the city by the year 2000, assuming the population continues to grow exponentially at a constant rate?The population of a particular city in 2000, assuming exponential growth at a constant rate, can be calculated based on the given information. The initial population in 1930 was 27,000, and the population in 1950 was 32,000. To find the growth rate, we can divide the population in 1950 by the population in 1930: 32,000 / 27,000 = 1.185185.
Now, using the formula for exponential growth, we can calculate the population in 2000. Let P(t) represent the population at time t, P(0) be the initial population, and r be the growth rate. The formula is P(t) = P(0) * [tex]e^(^r^t^)[/tex], where e is the mathematical constant approximately equal to 2.71828.
Plugging in the values, we have[tex]P(t) = 27,000 * e^(^1^.^1^8^5^1^8^5^*^7^0^)[/tex], where 70 represents the number of years from 1930 to 2000. Calculating this expression, we find P(t) ≈ 38,534.
Therefore, the population of the city in 2000 will be approximately 38,534 people.
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The population of the city in 2000 will be approximately 38,334 people.
To determine the population of the city in 2000, we can use the formula for exponential growth: P(t) = P₀ * e^(rt), where P(t) is the population at time t, P₀ is the initial population, e is the base of the natural logarithm (approximately 2.71828), r is the growth rate, and t is the time elapsed.
In this case, we have the initial population P₀ as 32,000 in 1950 and we need to find the population in 2000, which is a time span of 50 years. We can calculate the growth rate (r) using the formula: r = ln(P(t)/P₀) / t.
Plugging in the values, we have r = ln(38,334/32,000) / 50 ≈ 0.00825 (rounded to six decimal places). Now, substituting the known values into the exponential growth formula, we get P(2000) = 32,000 * e^(0.00825 * 50) ≈ 38,334 (rounded to the nearest whole number).
Therefore, the population of the city in 2000 will be approximately 38,334 people.
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If the median of a data set is 12 and the mean is 10, which of the following is most likely? Select the correct answer below:
O The data are skewed to the left.
O The data are skewed to the right.
O The data are symmetrical.
The median of a data set is 12 and the mean is 10, we need to determine the likely skewness of the data. The three options are: the data are skewed to the left, the data are skewed to the right, or the data are symmetrical.
When the median and the mean of a data set are not equal, it indicates that the data are skewed. Skewness refers to the asymmetry of the data distribution. If the median is greater than the mean, it suggests that the data are skewed to the left, also known as a left-skewed or negatively skewed distribution.
In this case, since the median is 12 and the mean is 10, the median is greater than the mean. This indicates that there is a tail on the left side of the distribution, pulling the mean towards lower values. Therefore, the data are most likely skewed to the left.
A left-skewed distribution typically has a long tail on the left side and a cluster of data points towards the right. This means that there are relatively more lower values in the data set compared to higher values.
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Given f(x, y) = 2y² + xy³ + 2ex, find f_xx
a. 2e^x
b. 3e^x
c. e^x
d. 6e^x
The solution to find f_xx is to take the second partial derivative of f(x, y) with respect to x, holding y constant. This gives f_xx = 2e^x.
To find f_xx, we first need to find the partial derivative of f(x, y) with respect to x. This gives f_x = y^3 + 2e^x.
Then, we take the partial derivative of f_x with respect to x. This gives f_xx = 2e^x.
Therefore, the answer is a. 2e^x.
Here is a more information of the steps involved:
1. We start by finding the partial derivative of f(x, y) with respect to x. This gives f_x = y^3 + 2e^x.
2. Then, we take the partial derivative of f_x with respect to x. This gives f_xx = 2e^x.
3. Therefore, the answer is a. 2e^x.
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A binomial distribution is composed of six fixed identical trials and the probability of success is 0,83. Therefore the mean and standard deviation of the binomial distribution are equal to: a. 4.98 and 0.9201 b. 1.02 and 0.9201
c. 1.50 and 0.866 d. 1.50 and 0.980
The binomial distribution that is composed of six identical fixed trials and a success probability of 0.83 has a mean and standard deviation of 4.98 and 0.9201, respectively. The correct option is A
The given probability distribution is a binomial distribution that consists of six identical fixed trials and the probability of success is 0.83.
Using the formula for the mean and standard deviation of the binomial distribution, we can solve this problem.
The formula for the mean and standard deviation is as follows:
Mean (μ) = [tex]n * p[/tex]
= [tex]6 * 0.83[/tex]
= 4.98
Standard deviation (σ) = √(n * p * q)
= √(6 * 0.83 * 0.17)
= 0.9201
Therefore, the mean and standard deviation of the binomial distribution are 4.98 and 0.9201, respectively. Thus, the correct option is (a)
The binomial distribution that is composed of six identical fixed trials and a success probability of 0.83 has a mean and standard deviation of 4.98 and 0.9201, respectively.
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i thought addition and subtraction can only be done from left to right (according to order of operations) but now they're grouping it? how do I solve this? what's the logic behind this? I'm confused:(
The two equivalent expressions are the ones at C and D.
-8/9 + 9/8
-(4/7 + 8/9) + 4/7 + 9/8
Which expressions are equivalent?Remember that for any sum, we have the associative property, which says that we can do a sum in any form:
A + B + C = A + (B + C) = (A + B) + C
So, here we have the sum:
-4/7 - 8/9 + 4/7 + 9/8
Using that property for the addition, we can group terms in any form we like, then the correct options are:
-(4/7 + 8/9) + 4/7 + 9/8
And we can also add the first term and the third ones, then we will get:
(-4/7 + 4/7) -8/9 + 9/8 = -8/9 + 9/8
Then the correct options are C and D.
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The one-to-one function is defined below. 6x f(x) = 4-5x Find f¹(x), where f¹ is the inverse of f. Also state the domain and range of f in interval notation.
The function f(x) = 4-5x is a one-to-one function. To find the inverse function f¹(x), we need to swap the roles of x and f(x) and solve for x.
To find the inverse function f¹(x), we swap the roles of x and f(x) in the equation f(x) = 4-5x. This gives us x = 4-5f¹(x). Solving this equation for f¹(x), we isolate f¹(x) to get f¹(x) = (4-x)/5.
The domain of f is the set of all possible values of x. In this case, there are no restrictions on x, so the domain is (-∞, +∞).
The range of f is the set of all possible values of f(x). By observing the equation f(x) = 4-5x, we see that f(x) can take any real number value. Therefore, the range is also (-∞, +∞) in interval notation.
In summary, the inverse function f¹(x) of f(x) = 4-5x is given by f¹(x) = (4-x)/5, and the domain and range of f are both (-∞, +∞).
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Suppose that x represents one of two positive numbers whose sum is 28. Determine a function f(x) that represents the product of these two numbers.
The function that would give the product of the numbers is f(x) = x (28 - x)
What is a function in mathematics?A function in mathematics is a relationship between a set of inputs (referred to as the domain) and a set of outputs (referred to as the codomain or range), where each input is connected to each output exactly once. Each input value is given a distinct output value.
We are told that the sum of the two numbers is 28 thus;
Let the first number be x
'Let the second number be 28 - x
We would have that;
f(x) = x (28 - x)
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Please kindly help with solving question
1. Find the exact value of each expression. Do not use a calculator. 5TT TT 7 TT 4 see (577) COS -√2 sin (177) 3 6 CSC
Evaluating the expression: 5TT TT 7 TT 4 see (577) COS -√2 sin (177) 3 6 CSC, the required exact value of the given expression is 2160° - 2√2 × sin (3°) + 1.
We know that TT = 180°. Hence, 5TT = 900°, 7TT = 1260°, and 4 see (577) = 4√3.
We know that cosine function is negative in the second quadrant, i.e., cos (θ) < 0 and sine function is positive in the third quadrant, i.e., sin (θ) > 0Hence, cos (177°) = -cos (180° - 3°) = -cos (3°) and sin (177°) = sin (180° - 3°) = sin (3°)
Using the trigonometric ratios of 30° - 60° - 90° triangle, we have CSC 30° = 2 and COT 30° = √3/3
Hence, COT 60° = 1/COT 30° = √3 and CSC 60° = 2 and TAN 60° = √3.
Now, we are ready to evaluate the expression.
5TT = 900°7TT = 1260°4 see (577) = 4√3cos (177°) = -cos (3°)sin (177°) = sin (3°)CSC 60° = 2COT 60° = √3CSC 30° = 2COT 30° = √3/3
∴ 5TT TT 7 TT 4 see (577) COS -√2 sin (177) 3 6 CSC = 900° + 1260° + 4√3 × (-1/√2) × sin (3°) + 3/6 × 2 = 2160° - 2√2 × sin (3°) + 1
The required exact value of the given expression is 2160° - 2√2 × sin (3°) + 1.
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7. Let a, b, c be integers, with a 0. Let ₁ and 2 be the roots of ax² + bx+c. (a) Show that if r₁ is rational, then so is 12. (b) Show that if a root is rational, then it can be written as, where p, q are integers, q divides a, and p divides c. (This is the Rational Roots Theorem for quadratic polynomials. You will need some facts from number theory to solve this problem.)
a) If r₁ is rational, then 12 is also rational.
b) If one of the roots is rational, then it can be written as p/q where p, q are integers, q divides a and p divides c.
Given that a, b, c are integers, with a ≠ 0. Let ₁ and 2 be the roots of
ax² + bx+c.
We need to show the following :
a) If r₁ is rational, then so is 12
b) If a root is rational, then it can be written as p/q where p, q are integers, q divides a and p divides c.
a) Let r₁ be rational.
Therefore, r₂= (b/a) - r₁ is also rational. Sum of roots ₁ and 2 is equal to -b/a.
Therefore,r₁ + r₂ = -b/a
=> r₂= -b/a - r₁
Now,
12= r₁ r₂
= r₁ (-b/a - r₁)
= -r₁² - (b/a) r₁
Therefore, if r₁ is rational, then 12 is also rational.
b) Let one of the roots be r.
Therefore,
ax² + bx+c
= a(x-r) (x-q)
= ax² - (a(r+q)) x + aqr
Now comparing the coefficients of x² and x, we get- (a(r+q))=b => r+q=-b/a ...(1) and
aqr=c
=> qr=c/a
=> q divides a and p divides c.
Now, substituting the value of q in equation (1), we get
r-b/a-q
=> r is rational.
Therefore, if one of the roots is rational, then it can be written as p/q where p, q are integers, q divides a and p divides c.
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You don't need problem 6. It just needs the answer to be in a piecewise function. Sorry for the confusion.
Let x = 100+ 100fe. Plot y = x-100? 100£ over the interval 0 ≤ f≤ 1.
a) Describe the result as a piecewise function as in P6.
b) Explain (XC).
(c) What is the advantage of this method of computing £?
The result can be described as a piecewise function:
```
y = 0, if 0 ≤ f < 0.01
y = 100, if 0.01 ≤ f ≤ 1
```
What does (XC) refer to in the context of this problem?The advantage of using a piecewise function to compute £ is that it allows for different calculations based on the value of the variable f. By defining different cases for the function, we can handle specific ranges of f differently, resulting in a more accurate and flexible computation. This method allows us to assign a constant value to y within each range, simplifying the calculations and providing a clear representation of the relationship between x and y. It helps to capture the behavior of the function over the given interval and provides a structured approach to handling different scenarios.
y = 0, if 0 ≤ f < 0.01
y = 100, if 0.01 ≤ f ≤ 1
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The variable ‘AgencyEngagement’ is a scale measurement that indicates how engaged an employee is with their Agency/Department. This variable was measured on a scale that can take values from 0 to 20, with higher values representing greater employee engagement with their Agency/Department. Produce the relevant graph and tables to summarise the AgencyEngagement variable and write a paragraph explaining the key features of the data observed in the output in the style presented in the course materials. Produce the relevant graph and tables to summarise the ‘AgencyEngagement’ variable and write a paragraph explaining the key features of the data observed in the output in the style presented in the course materials. Which is the most appropriate measure to use of central tendency, that being node median and mean?
To summarize the 'AgencyEngagement' variable, we can create a graph and tables. Additionally, we need to determine whether it is the mode, median, or mean.
To summarize the 'AgencyEngagement' variable, we can start by creating a histogram or bar graph that shows the frequency or count of each engagement score on the x-axis and the number of employees on the y-axis. This graph will provide an overview of the distribution of engagement scores and any patterns or trends in the data.
Additionally, we can create a table that displays summary statistics for the 'AgencyEngagement' variable. This table should include measures of central tendency (mean, median, and mode), measures of dispersion (range, standard deviation), and any other relevant statistics such as minimum and maximum values.
In analyzing the key features of the data observed in the output, we should examine the shape of the distribution. If the distribution is approximately symmetric, then the mean would be an appropriate measure of central tendency. However, if the distribution is skewed or contains outliers, the median may be a better measure as it is less influenced by extreme values. The mode can also provide insights into the most common level of engagement.
Therefore, to determine the most appropriate measure of central tendency for the 'AgencyEngagement' variable, we need to assess the shape of the distribution and consider the presence of outliers. If the distribution is roughly symmetrical without significant outliers, the mean would be suitable. If the distribution is skewed or has outliers, the median should be used as it is more robust. Additionally, the mode can provide information about the most prevalent level of engagement.
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6. An airplane is headed north with a constant velocity of 430 km/h. the plane encounters a west wind blowing at 100 km/h. a) How far will the plane travel in 2 h? b) What is the direction of the plan
The direction of the plane is still north, because the plane is moving forward at a greater speed than the wind is pushing it back.
a) The plane will travel 760 km in 2 hours. To solve this, we need to first calculate the resultant velocity of the plane.
The resultant velocity is 430 km/h in the northwards direction plus the wind velocity of 100 km/h in the westwards direction.
This results in a velocity vector of $(430)² + (100)² = 468.3$ km/h in the northwest direction.
As the plane has a velocity of 468.3 km/h in this direction, it will travel $(468.3)(2)$ = 936.6 km in 2 hours.
b) The direction of the plane is northwest.
Therefore, the direction of the plane is still north, because the plane is moving forward at a greater speed than the wind is pushing it back.
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Let U = {a, b, c, d, e, f, g, h, i, j, k}, A = {a, f, g, h, j, k}, B = {a, b, g, h, k} C = {b, c, f, j, k} Determine AU ( CB). Select the correct choice and, if necessary, fill in the answer box to complete your choice. O A. AU (COB)' = (Use a comma to separate answers as needed.) OB. AU (COB) is the empty set.
The AU (CB)' = U - AU (CB) = {c, d, e, i}We can see that option A, AU (CB)' = {c, d, e, i}, is the correct answer.The union of two sets A and B, denoted by A ∪ B
Let U = {a, b, c, d, e, f, g, h, i, j, k}, A = {a, f, g, h, j, k}, B = {a, b, g, h, k} C = {b, c, f, j, k}. We need to determine AU ( CB).Solution:
, is the set that contains those elements that are either in A or in B or in both.
That is,A ∪ B = {x : x ∈ A or x ∈ B}The intersection of two sets A and B, denoted by A ∩ B, is the set that contains those elements that are in both A and B.
That is,A ∩ B = {x : x ∈ A and x ∈ B}AU (CB) = {x : x ∈ A or x ∈ (C ∩ B)} = {a, f, g, h, j, k} ∪ {b, k} = {a, b, f, g, h, j, k}CB = {x : x ∈ C and x ∈ B} = {g, h, k}
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Suppose that a game has a payoff matrix
A = [\begin{array}{cccc}-20&30&-20&1\\21&-31&11&40\\-40&0&30&-10\end{array}\right]
If players R and C use strategies
p = [\begin{array}{ccc}1/2&0&1/2\end{array}\right] and
q = [\begin{array}{c}1/4\\1/4\\1/4\end{array}\right]
respectively, what is the expected payoff of the game? E(p, q) =
The expected payoff of the game with strategies p and q is 1.875.To calculate the expected payoff of the game with the given strategies, we need to multiply the payoff matrix A with the strategy vectors p and q.
Let's perform the matrix multiplication:
A * p = [\begin{array}{cccc}-20&30&-20&1\\21&-31&11&40\\-40&0&30&-10\end{array}\right] * [\begin{array}{ccc}1/2\\0\\1/2\end{array}\right]
= [\begin{array}{c}-20*(1/2) + 30*(0) - 20*(1/2) + 1*(1/2)\\21*(1/2) - 31*(0) + 11*(1/2) + 40*(1/2)\\-40*(1/2) + 0*(0) + 30*(1/2) - 10*(1/2)\end{array}\right]
= [\begin{array}{c}-10 + 0 - 10 + 1/2\\10.5 + 0 + 5.5 + 20\\-20 + 0 + 15 - 5\end{array}\right]
= [\begin{array}{c}-18.5\\36\\-10\end{array}\right]
Now, let's calculate the dot product of the result with the strategy vector q:
E(p, q) = [\begin{array}{ccc}-18.5&36&-10\end{array}\right] * [\begin{array}{c}1/4\\1/4\\1/4\end{array}\right]
= -18.5*(1/4) + 36*(1/4) - 10*(1/4)
= -4.625 + 9 - 2.5
= 1.875
Therefore, the expected payoff of the game with strategies p and q is 1.875.
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77. Find the inverse of the nonsingular matrix -4 1 6 -2]
The inverse of the nonsingular matrix [-4 1; 6 -2] is [1/2 1/2; -3/4 -1/4].
To find the inverse of a matrix, we follow a specific procedure. Let's consider the given matrix [-4 1; 6 -2] and find its inverse.
Step 1: Calculate the determinant of the matrix.
The determinant of the matrix is found by multiplying the diagonal elements and subtracting the product of the off-diagonal elements. For the given matrix, the determinant is:
Det([-4 1; 6 -2]) = (-4) * (-2) - (1) * (6) = 8 - 6 = 2.
Step 2: Determine the adjugate matrix.
The adjugate matrix is obtained by taking the transpose of the matrix of cofactors. To find the cofactors, we interchange the signs of the elements and compute the determinants of the remaining 2x2 matrices. For the given matrix, the cofactor matrix is:
[-2 -6; -1 -4].
Taking the transpose of this matrix, we get the adjugate matrix:
[-2 -1; -6 -4].
Step 3: Calculate the inverse matrix.
The inverse of the matrix is obtained by dividing the adjugate matrix by the determinant. For the given matrix, the inverse is:
[1/2 1/2; -3/4 -1/4].
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Consider the problem maxx +2y subject to x² + y² ≤ 1 and x + y ≥ 0 a. Write down the first order conditions. b. Solve the problem.
The problem involves maximizing the objective function f(x, y) = x + 2y, subject to the constraints x² + y² ≤ 1 and x + y ≥ 0.
In order to solve the problem, we need to determine the first-order conditions and find the optimal solution.
a. First-order conditions:
To find the first-order conditions, we need to consider the Lagrangian function L(x, y, λ) = f(x, y) - λ(g(x, y)), where g(x, y) represents the constraints. In this case, the constraints are x² + y² ≤ 1 and x + y ≥ 0.
The first-order conditions are:
∂L/∂x = 1 - 2λx = 0
∂L/∂y = 2 - 2λy = 0
g(x, y) = x² + y² - 1 ≤ 0
h(x, y) = -(x + y) ≤ 0
b. Solving the problem:
To solve the problem, we need to solve the first-order conditions and check the feasibility of the constraints.
From the first-order conditions, we have:
1 - 2λx = 0 --> x = 1/(2λ)
2 - 2λy = 0 --> y = 1/(2λ)
Substituting these values into the constraint equations, we have:
(1/(2λ))² + (1/(2λ))² ≤ 1 --> 1/(4λ²) + 1/(4λ²) ≤ 1 --> 1/λ² ≤ 1 --> λ² ≥ 1 --> λ ≥ 1 or λ ≤ -1
Since λ must be non-negative, we have λ ≥ 1.
Substituting λ = 1 into the expressions for x and y, we get:
x = 1/2
y = 1/2
Therefore, the optimal solution is x = 1/2 and y = 1/2, which maximizes the objective function x + 2y subject to the given constraints.
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1- Find the domain of the function. (Enter your answer using interval notation.) H(t) = 81 − t2/ 9 − t. Sketch graph of the function.
2- Find the domain of the function. (Enter your answer using interval notation.) Sketch a graph of this fuction.
f(x) =
3 −
1
2
x if x ≤ 2
9x − 2 if x > 2
3- Sketch the graph of the function.
f(x) =
To find the domain of the function H(t) = (81 - t^2) / (9 - t), we need to consider the values of t that make the denominator (9 - t) non-zero since division by zero is undefined.
First, let's find the values that make the denominator zero:
9 - t = 0
t = 9
So, t = 9 is not in the domain of the function H(t) because it would result in division by zero.
Therefore, the domain of the function H(t) is (-∞, 9) U (9, +∞).
To sketch the graph of the function H(t), we start by plotting some key points on the graph. Here are a few points you can plot:
Choose some values for t in the domain, such as t = -10, -5, 0, 5, 8, and 10.
Calculate the corresponding values of H(t) using the given function.
Plot the points (-10, H(-10)), (-5, H(-5)), (0, H(0)), (5, H(5)), (8, H(8)), and (10, H(10)).
Connect the plotted points smoothly to form the graph. Keep in mind that the graph will have an asymptote at t = 9 because of the denominator being zero at that point.
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A patient needs 3 L of D5W with 20 meq of potassium chloride to infuse over one day (24 hours). The DF is 15 gtt/mL. What is the correct rate of flow in gtt/min? Round to the nearest whole number.
The correct rate of flow in gtt/min for infusing 3 L of D5W with 20 meq of potassium chloride over 24 hours is 31 gtt/min.
To determine the rate of flow in gtt/min, we need to calculate the total number of drops needed over the infusion period and then divide it by the total time in minutes.
First, we need to find the total volume of the solution in milliliters (mL):
3 L = 3000 mL
Next, we calculate the total number of drops needed. We can use the drop factor (DF) of 15 gtt/mL:
Total drops = Volume (mL) x DF
Total drops = 3000 mL x 15 gtt/mL
Next, we calculate the total time in minutes:
24 hours = 24 x 60 minutes = 1440 minutes
Finally, we divide the total drops by the total time in minutes to find the rate of flow in gtt/min:
Rate of flow (gtt/min) = Total drops / Total time (minutes)
Rate of flow (gtt/min) = (3000 mL x 15 gtt/mL) / 1440 minutes
Simplifying the expression, we have:
Rate of flow (gtt/min) ≈ 31.25 gtt/min
Rounding to the nearest whole number, the correct rate of flow in gtt/min is approximately 31 gtt/min.
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С x 4 Gx 2 + y2 = Being the curre from Point (30) to point Co-3) on the circle 9 as sy 2 ds c Calculate the Integral)
The solution of the integral is ∫[C] (c * 4 * G * 2 + y * 2) ds = ∫[arcsin(-1/3), arccos(1/3)] (c * 4 * G * 2 + 9 * sin(t)²) * 9 dt
To calculate the integral of the given expression over the curve on the circle, we first need to parameterize the curve. Let's denote the parameter along the curve as t. We can represent the curve on the circle as (x(t), y(t)), where x(t) and y(t) are the x-coordinate and y-coordinate of the curve at parameter t.
Since the curve lies on the circle with center C and radius 9, we can use the equation of a circle to find x(t) and y(t). The equation of a circle with center (a,b) and radius r is given by:
(x - a)² + (y - b)² = r²
In our case, the center C is (0,0) and the radius is 9. Plugging in these values, we have:
x(t)² + y(t)² = 9²
Next, let's solve for x(t) and y(t) in terms of the parameter t. One way to parameterize the curve on the circle is by using trigonometric functions. We can express x(t) and y(t) as:
x(t) = 9 * cos(t) y(t) = 9 * sin(t)
Now that we have the parameterization of the curve, we can calculate the line integral. The line integral of a function f(x, y) over a curve C parameterized by x(t) and y(t) is given by:
∫[C] f(x, y) ds = ∫[a,b] f(x(t), y(t)) * ||r'(t)|| dt
In this case, the function we want to integrate is c * 4 * G * 2 + y * 2, where c and G are constants. Plugging in the parameterization of the curve, we have:
∫[C] (c * 4 * G * 2 + y * 2) ds = ∫[a,b] (c * 4 * G * 2 + 9 * sin(t)²) * ||r'(t)|| dt
To calculate ||r'(t)||, we differentiate x(t) and y(t) with respect to t:
x'(t) = -9 * sin(t) y'(t) = 9 * cos(t)
The magnitude of the derivative vector r'(t) is given by ||r'(t)|| = √(x'(t)² + y'(t)²). Plugging in the values, we have:
||r'(t)|| = √((-9 * sin(t))² + (9 * cos(t))²) = √(81 * sin(t)² + 81 * cos(t)²) = √(81) = 9
Therefore, the line integral simplifies to:
∫[C] (c * 4 * G * 2 + y * 2) ds = ∫[a,b] (c * 4 * G * 2 + 9 * sin(t)²) * 9 dt
Now, we need to determine the limits of integration. We are given that the curve starts at point (3,0) and ends at point (0,-3). We can find the values of t that correspond to these points by plugging the values of x and y into the parameterization equations:
When x = 3 and y = 0: 3 = 9 * cos(t) => cos(t) = 1/3 => t = arccos(1/3)
When x = 0 and y = -3: -3 = 9 * sin(t) => sin(t) = -1/3 => t = arcsin(-1/3)
Therefore, the limits of integration are a = arcsin(-1/3) and b = arccos(1/3).
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(a) Let f(x) = x² + 1. Compute f(0), ƒ(-1), f(1/2), and f(√2).
(b) For what values of x is it true that
(i) f(x) = f(-x)?
(ii) f(x + 1) = f(x) + f(1)?
(iii) f(2x) = 2ƒf(x)?
Problem 2
The cost of producing x units of a commodity is given by C(x) = 1000 + 300x + x².
(a) Compute C(0), C(100), and C(101) - C(100).
(b) Compute C(x + 1) - C(x), and explain in words the meaning of the difference.
For problem 1,
we are given f(x) = x² + 1.
The
values
of f(0), f(-1), f(1/2), and f(√2) are 1, 2, 1.25, and 3, respectively.
For problem 2,
We are given C(x) = 1000 + 300x + x².
The
marginal cost
is constant at 300.
We are given f(x) = x² + 1
Let’s compute the values of x for which the following hold true:
(i) f(x) = f(-x)
x² + 1 = (-x)² + 1 x²
=x²
Therefore, the above holds true for all x.
(ii) f(x + 1) = f(x) + f(1) (x + 1)² + 1
=x² + 1 + 1² + 1 x² + 2x + 1 + 1
= x² + 2 2x
= 0 x
= 0
Therefore, the above holds true only for x = 0.
(iii) f(2x) = 2f(x) (2x)² + 1
= 2(x² + 1) 4x² + 1
= 2x² + 2 2x²
= 1 x
= ± 1/√2
Therefore, the above holds true for x = 1/√2 and
x = -1/√2
(i) f(x) = f(-x) holds
true
for all x.
(ii) f(x + 1)
= f(x) + f(1) holds true only for
x = 0.
(iii) f(2x) = 2f(x) holds true for
x = 1/√2 and
x = -1/√2.
We are given C(x) = 1000 + 300x + x².
C(x + 1) – C(x) = [1000 + 300(x + 1) + (x + 1)²] – [1000 + 300x + x²] C(x + 1) – C(x)
= 300 + 2x
The above difference gives the marginal cost of producing one extra unit of the
commodity
.
The marginal cost is a constant value of 300, whereas, 2x is the variable cost associated with the
production
of an additional unit of the commodity.
C(x + 1) – C(x) gives the marginal cost of producing one extra unit of the commodity.
The marginal cost is constant at 300, whereas 2x is the variable cost associated with the production of an additional unit of the commodity.
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Let H be a Hilbert space. From Riesz' theorem we know that the conjugate linear map
L: H→H', v (ov: w→ (v, w))
is an isometry.
(a) Use this map L to find a canonical conjugate linear isometry K: H'H".
(b) Show that KoL=j: H→ H", the canonical inclusion into the bidual space defined by j(x): o→ o(x).
The canonical conjugate linear isometry K: H'H" can be obtained by composing the conjugate linear map L: H→H' with the canonical conjugate linear map J: H'→H". The resulting map K is an isometry. The equality KoL = j holds, where j is the canonical inclusion map from H to H", as J(L(v)) = L(v) = v'' for any element v in H.
a) To compute the canonical conjugate linear isometry K: H'H", we can compose the conjugate linear map L: H→H' with the canonical conjugate linear map J: H'→H". The composition K = J∘L gives us the desired map K: H'H" defined by K(v')(w'') = L(v')(J(w'')). This map K is an isometry.
(b) To show that KoL = j: H→H", we need to demonstrate that for any element v in H, the image of v under KoL is equal to the image of v under j.
Using the definition of K from part (a), we have KoL(v) = K(L(v)) = J(L(v)). On the other hand, the image of v under j is j(v) = v''.
To establish the equality KoL = j, we need to show that J(L(v)) = v''. Since J is the canonical inclusion map from H' to H", it maps elements of H' to their corresponding elements in H".
Since L(v) is an element of H', we can identify J(L(v)) with L(v) in H". Therefore, J(L(v)) = L(v) = v''.
Thus, we have shown that KoL = j, confirming the equality between the composition of the maps K and L and the canonical inclusion map j.
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