Prove that if S and T are isomorphic submodules of a module M it does not necessarily follow that the quotient modules M/S and M/T are isomorphic. Prove also that if ST ST₂ as modules it does not necessarily follow that T₁ T₂. Prove that these statements do hold if all modules are free and have finite rank.

Given differential equation is:2"x(t) - 42"x(t) + 4r(t) = 0Given differential equation is a second order linear homogeneous coordinates  

differential equation, whose characteristic equation is:2m² - 42m + 4 = 0 ⇒ m² - 21m + 2 = 0Solving above quadratic equation, we get:m₁ = 20.9282 and m₂ = 0.0718So,

the general solution of the given differential equation can be written as:x(t) = C₁e⁽²⁰.⁹²⁸²t⁾ + C₂e⁽⁰.⁰⁷¹⁸t⁾Where C₁ and C₂ are constants of integration.To find the solution of the differential

This is the answer to the given problem.

This answer is a as we have to solve the given differential equation using the standard method of finding the general solution of second order linear homogeneous differential equation and then find the solution of the differential equation associated with the given initial conditions.

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a) No, b) Yes, c) Yes, d) No, e) No, f) Yes, g) Yes, h) Yes, i) Yes, j) Yes. In (AB)∗ is a concatenation of zero or more strings from AB, which is exactly the definition of A∗B∗.

a) The statement A(BC) ⊆ (AB)C is not always true. A counterexample is when A = {a}, B = {b}, and C = {c}. In this case, A(BC) = {abc}, while (AB)C = {(ab)c} = {abc}. Therefore, A(BC) = (AB)C, and the statement is false.

b) The statement A(BC) ⊇ (AB)C is always true. This is because the left-hand side contains all possible concatenations of a string from A, a string from B, and a string from C, while the right-hand side contains only the concatenations where the string from A is concatenated with the concatenation of strings from B and C.

c) The statement A(B ∪ C) ⊆ AB ∪ AC is always true. This is because any string in A(B ∪ C) is a concatenation of a string from A and a string from either B or C, which is exactly the definition of AB ∪ AC.

d) The statement A(B ∪ C) ⊇ AB ∪ AC is not always true. A counterexample is when A = {a}, B = {b}, and C = {c}. In this case, A(B ∪ C) = A({b, c}) = {ab, ac}, while AB ∪ AC = {ab} ∪ {ac} = {ab, ac}. Therefore, A(B ∪ C) = AB ∪ AC, and the statement is false.

e) The statement A(B ∩ C) ⊆ AB ∩ AC is not always true. A counterexample is when A = {a}, B = {b}, and C = {c}. In this case, A(B ∩ C) = A({}) = {}, while AB ∩ AC = {ab} ∩ {ac} = {}. Therefore, A(B ∩ C) = AB ∩ AC, and the statement is false.

f) The statement A(B ∩ C) ⊇ AB ∩ AC is always true. This is because any string in AB ∩ AC is a concatenation of a string from A and a string from both B and C, which is exactly the definition of A(B ∩ C).

g) The statement A∗ ∪ B∗ ⊆ (A ∪ B)∗ is always true. This is because A∗ ∪ B∗ contains all possible concatenations of zero or more strings from A or B, while (A ∪ B)∗ also contains all possible concatenations of zero or more strings from A or B.

h) The statement A∗ ∪ B∗ ⊇ (A ∪ B)∗ is always true. This is because any string in (A ∪ B)∗ is a concatenation of zero or more strings from A or B, which is exactly the definition of A∗ ∪ B∗.

i) The statement A∗B∗ ⊆ (AB)∗ is always true. This is because A∗B∗ contains all possible concatenations of zero or more strings from A followed by zero or more strings from B, while (AB)∗ also contains all possible concatenations of zero or more strings from AB.

j) The statement A∗B∗ ⊇ (AB)∗ is always true. This is because any string

in (AB)∗ is a concatenation of zero or more strings from AB, which is exactly the definition of A∗B∗.

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The sum of the simple probabilities for a collectively exhaustive set of outcomes must be equal to one, serving as a fundamental principle of probability theory. This principle holds true for any situation where events are mutually exclusive and cover all possible outcomes.

The sum of the simple probabilities for a collectively exhaustive set of outcomes must be equal to one.

This fundamental principle is a cornerstone of probability theory and ensures that all possible outcomes are accounted for.

To understand why the sum of probabilities must equal one, let's consider a simple example. Imagine flipping a fair coin.

The two possible outcomes are "heads" and "tails." Since these two outcomes cover all possibilities, they form a collectively exhaustive set. The probability of getting heads is 0.5, and the probability of getting tails is also 0.5.

When we add these probabilities together (0.5 + 0.5), we get 1, indicating that the sum of probabilities for the complete set of outcomes is indeed one.

This principle extends beyond coin flips to any situation involving mutually exclusive and collectively exhaustive events.

For instance, if we roll a standard six-sided die, the probabilities of getting each face (1, 2, 3, 4, 5, or 6) are all 1/6.

When we add these probabilities together (1/6 + 1/6 + 1/6 + 1/6 + 1/6 + 1/6), we again obtain 1.

The requirement for the sum of probabilities to equal one ensures that the total probability space is accounted for, leaving no room for events outside of it.

It provides a mathematical framework for reasoning about uncertain events and allows us to quantify the likelihood of various outcomes.

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The letter "F" is used to refer to the theory-based standardized statistic for comparing several means. So, correct option is D.

The F-statistic is commonly used in statistical analysis to determine whether the means of two or more groups are significantly different from each other.

The F-statistic is derived from the F-distribution, which is a probability distribution that arises when comparing variances or ratios of variances. In the context of comparing means, the F-statistic is calculated by dividing the variance between groups by the variance within groups.

By comparing the calculated F-statistic to critical values from the F-distribution, we can determine whether there is a significant difference between the means of the groups being compared. If the calculated F-statistic is larger than the critical value, it suggests that there is a significant difference between at least two of the means.

Therefore, when comparing several means and conducting hypothesis tests or analysis of variance (ANOVA), the letter "F" is used to represent the theory-based standardized statistic.

So, correct option is D.

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The inverse of matrix A is:

[tex][\left[\begin{matrix}1.5&2.5&-1\\-2.5&-4.5&2\\-0.5&-0.5&1\end{matrix}\right]\][/tex]

The augmented matrix of the system of equations is:

[tex]| 2 -3 0 1 || 1 0 -1 0 || 1 1 1 5 |[/tex]

Now, we are going to use elementary row operations to solve this system of equations.

First, let's multiply R1 by 1/2 to get a leading 1 in R1.

[tex]| 1 -3/2 0 1/2 || 1 0 -1 0 || 1 1 1 5 |[/tex]

Next, we want to use R1 to get zeros under the leading 1 in R1.

[tex]| 1 -3/2 0 1/2 || 0 3/2 -1/2 -1/2 || 0 3/2 1/2 9/2 |[/tex]

Now, we want to use elementary row operations to get zeros in the third row of the matrix.

[tex]| 1 -3/2 0 1/2 || 0 3/2 -1/2 -1/2 || 0 0 1 5 |[/tex]

We will back substitute to get values for y and x.

[tex]| 1 -3/2 0 1/2 || 0 1 0 2 || 0 0 1 5 |x = -2y + 1z = 5[/tex]

Now, let's write the system of equations in Aữ = b form:[tex]2x - 3y + 0z = 1x + 0y - z = 0x + y + z = 5\[A\] =  \[\left[\begin{matrix}2&-3&0\\1&0&-1\\1&1&1\end{matrix}\right]\]\[u\] =  \[\left[\begin{matrix}x\\y\\z\end{matrix}\right]\]\[b\] = \[\left[\begin{matrix}1\\0\\5\end{matrix}\right]\][/tex]

Find the inverse of matrix A from the question.

[tex]| 2 -3 0 || 1 0 -1 || 1 1 1 |[/tex]

Now, we will use elementary row operations to get an identity matrix on the left side of the matrix.

[tex]| 1 0 0 || 13/2 1 0 || 3/2 5 -2 || -5/2 0 1 |[/tex]

The inverse of matrix A is:

[tex][\left[\begin{matrix}1.5&2.5&-1\\-2.5&-4.5&2\\-0.5&-0.5&1\end{matrix}\right]\][/tex]

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Any value of b that is equal to or less than 0.5 (half the total volume) would satisfy the condition.  The glass is half full: 50% volume.

"The glass is half full" is a metaphorical expression used to describe an optimistic or positive perspective. It suggests focusing on the portion of a situation that is favorable or has been accomplished, rather than dwelling on what is lacking or incomplete.

In this exercise, if the glass is filled to a height of b = 7 cm, we need to calculate the percentage of the total volume this represents. To do so, we compare the volume for 7 cm (V7) with the volume for 14 cm (V14) and express it as a percentage.

The volume of the glass filled to a height of b = 7 cm is half the volume when it is filled to the top, which means V7 = 0.5 * V14.

To find the percentage, we can use the formula (V7 / V14) * 100

By substituting V7 = 0.5 * V14 into the formula, we have (0.5 * V14 / V14) * 100 = 0.5 * 100 = 50%.

Therefore, if the glass is filled to a height of b = 7 cm, it represents 50% of the total volume.

Now, let's calculate the volume contained in the glass when it is filled to the top, b = 14 cm. The volume is given as 1, in the exercise.

To convert the volume from cm³ to ounces, we divide by 1000 and multiply by 33.82. So, the volume in ounces would be (1 / 1000) * 33.82 = 0.03382 ounces.

Finally, to find a value of b within 1 decimal point that gives half the volume when the glass is full, we can set up the equation Vb = 0.5 * V14 and solve for b.

0.5 * V14 = 1 * V14

0.5 = V14

Therefore, any value of b that is equal to or less than 0.5 (half the total volume) would satisfy the condition.

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The hypothesis test aims to determine whether female college students score higher than 3.0 on the emotional empathy scale. The null hypothesis states that there is no significant difference, while the alternative hypothesis suggests that there is a significant difference.

a. The null hypothesis (H₀) states that the mean emotional empathy score for female college students is equal to or less than 3.0 (μ ≤ 3.0), while the alternative hypothesis (H₁) proposes that the mean emotional empathy score for female college students is greater than 3.0 (μ > 3.0). To compute the test statistic, we use the formula:

t = (sample mean - population mean) / (population standard deviation / √sample size)

In this case, the sample mean response is 3.28, the population mean is 3.0, the population standard deviation is 0.5, and the sample size is 30. Plugging these values into the formula, we calculate the test statistic. To find the observed significance level (p-value) of the test, we compare the test statistic to the appropriate t-distribution with (sample size - 1) degrees of freedom. By looking up the p-value associated with the test statistic in the t-distribution table or using statistical software, we determine the significance level.

With a significance level of α = 0.01, we compare the observed significance level (p-value) from part c to α. If the p-value is less than α, we reject the null hypothesis. Otherwise, we fail to reject the null hypothesis. The choice of significance level α depends on the desired level of confidence in the results. The smaller the α-value, the stronger the evidence required to reject the null hypothesis. As long as the observed significance level (p-value) is smaller than the chosen α-value, we can reject the null hypothesis and conclude that there is sufficient evidence to support the alternative hypothesis.

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The answer is 0.7.The calculation of PA or B) has been presented above, and it is equal to 0.7.

PA and B represents the intersection of A and B, meaning the probability of A and B happening simultaneously. PA or B means the union of A and B, i.e., the probability of A or B happening.

The following formula can be used to calculate it: P(A or B) = P(A) + P(B) - P(A and B)Using the given values, we can substitute them into the formula to calculate the probability of A or B happening:P(A or B) = P(A) + P(B) - P(A and B)P(A or B) = 0.3 + 0.6 - 0.2P(A or B) = 0.7The probability of A or B happening is 0.7.

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a-1. The value of the population mean is unknown.a-2. The best estimate of this value is 24c. The value of t for a 90% confidence level can be calculated using the t-distribution table. Since the sample size is less than 30 and the population standard deviation is unknown, a t-distribution is used.

Using a t-distribution table with 18 degrees of freedom (n - 1)

The value of t for a 90% confidence level is 1.734 (approx.).

d. The margin of Error is calculated as follows:

M.E. = t * (s/√n)

Where, t = 1.734 (from part c)

s = 4 (standard deviation)

n = 19 (sample size)

M.E. = 1.734 * (4/√19)M.E. = 1.734 * 0.918M.E. = 1.59012 ≈ 1.59

Therefore, the margin of error is 1.59

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The Taylor series is [tex]ln(x) = x - x^2/2 + x^3/3 - x^4/4 + ...[/tex] , The Fourier series is  [tex]f(t) = (1 - cos(t))/2 + 9/(2\pi) sin(t)[/tex] , The transfer function is[tex]H(s) = (35s-140)/((5s+2)(s-5))[/tex].

The Taylor series for the function[tex]f(x) = ln(x)[/tex] about x = 0 can be found using the following steps:

Let [tex]f(x) = ln(x)[/tex].

Let [tex]f(0) = ln(1) = 0[/tex].

Let[tex]f'(x) = 1/x[/tex].

Let[tex]f''(x) = -1/x^2[/tex].

Continue differentiating f(x) to find higher-order derivatives.

Substitute x = 0 into the Taylor series formula to get the Taylor series for f(x) about x = 0.

The Taylor series for[tex]f(x) = ln(x)[/tex] about x = 0 is:

[tex]ln(x) = x - x^2/2 + x^3/3 - x^4/4 + ...[/tex]

The Fourier series of the function [tex]f(t) = {-1; -\pi \leq t \leq 0 0 < t < \pi 19 1}[/tex]can be found using the following steps:

Let [tex]f(t) = {-1; -\pi \leq t \leq 0 0 < t < \pi 19 1}[/tex].

Let [tex]a_0 = 1/2[/tex].

Let[tex]a_1 = -1/(2\pi)[/tex].

Let [tex]a_2 = 9/(2\pi^2).[/tex]

Let[tex]b_0 = 0[/tex].

Let[tex]b_1 = 1/(2\pi)[/tex].

Let[tex]b_2 = 0.[/tex]

The Fourier series for f(t) is:

[tex]f(t) = a_0 + a_1cos(t) + a_2cos(2t) + b_1sin(t) + b_2sin(2t)[/tex]

[tex]= (1 - cos(t))/2 + 9/(2\pi) sin(t)[/tex]

The transfer function[tex]H(s) = 7/(5s+2) + 3/(s-5)[/tex]can be found using the following steps:

Let [tex]H(s) = 7/(5s+2) + 3/(s-5).[/tex]

Find the partial fraction decomposition of H(s).

The transfer function is the ratio of the numerator polynomial to the denominator polynomial.

The partial fraction decomposition of [tex]H(s) = 7/(5s+2) + 3/(s-5)[/tex] is:

[tex]H(s) = (7/(5(s-5))) + (3/(s-5))\\= (7/5) (1/(s-5)) + (3/5) (1/(s-5))\\= (2) (1/(s-5))[/tex]

The transfer function is:

[tex]H(s) = (2)/(s-5)[/tex]

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The first three members of the corresponding orthonormal basis of V are:

[tex]v0(x) = 1, \\v1(x) = sqrt(2) x, \\v2(x) = 2x2 - 1.[/tex]

Given: V be the Euclidean space of polynomials with the inner product [tex](u, v) S* w(x)u(x)v(x)dx[/tex] where [tex]w(x) = xe-r[/tex].

With [tex]Un(x) = x", \\n = 0, 1, 2, ...[/tex]

To determine: the first three members of the corresponding orthonormal basis of VFormula to find

Orthonormal basis of V is: {vi}, where for each [tex]= sqrt((ui,ui)).i.e {vi} = {ui(x)/sqrt((ui,ui))}[/tex]

with ||ui|| [tex]= sqrt((ui,ui)).i.e {vi} \\= {ui(x)/sqrt((ui,ui))}[/tex]

, where ([tex]ui,uj) = S*w(x)ui(x)uj(x)dx[/tex]

Here w(x) = xe-r and Un(x) = xn

First we find the inner product of U[tex]0(x), U1(x) and U2(x).\\S* w(x)U0(x)U0(x)dx = S* xe-r (1)(1)dx=[/tex]

integral from 0 to infinity (xe-r dx)= x (-e-r x - 1) from 0 to infinity

[tex]= 1S* w(x)U1(x)U1(x)dx \\= S* xe-r (x)(x)dx=[/tex]

integral from 0 to infinity

[tex](x2e-r dx)= 2S* w(x)U2(x)U2(x)dx \\= S* xe-r (x2)(x2)dx=[/tex]

integral from 0 to infinity[tex](x4e-r dx)= 24[/tex]

We have

[tex](U0,U0) = 1, \\(U1,U1) = 2, \\(U2,U2) = 24[/tex]

So the corresponding orthonormal basis of V are:

[tex]v0(x) = U0(x)/||U0(x)|| = 1, \\v1(x) = U1(x)/||U1(x)|| = sqrt(2) x, \\v2(x) = U2(x)/||U2(x)|| \\= sqrt(24/6) (x2 - (1/2))\\= sqrt(4) (x2 - (1/2))\\= 2x2 - 1[/tex]

Therefore, the first three members of the corresponding orthonormal basis of V are

[tex]v0(x) = 1, \\v1(x) = sqrt(2) x, \\v2(x) = 2x2 - 1.[/tex]

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The substitution v = p(x)y', where p(x) is a suitable function, the self-adjoint second-order differential equation can be reduced to the special Riccati equation.

To demonstrate the reduction of the self-adjoint second-order differential equation into the special Riccati equation, we begin with the given equation:

(d/dx)(p(x)y') + q(x)y = 0

First, we differentiate v = p(x)y' with respect to x:

dv/dx = d/dx(p(x)y')

Using the product rule, we can expand the derivative:

dv/dx = p'(x)y' + p(x)y''

Now, substituting v = p(x)y' into the original equation, we have:

(dv/dx) + q(x)y = p'(x)y' + p(x)y'' + q(x)y = 0

Rearranging the terms, we obtain:

p(x)y'' + (p'(x) + q(x))y' + q(x)y = 0

Comparing this equation with the general form of the Riccati equation:

[tex](du/dx) + a(x)u + b(x)u^2 + c(x) = 0[/tex]

We can identify the coefficients as follows:

[tex]a(x) = (p'(x) + q(x))/p(x)b(x) = 0 (no u^2 term in the reduced equation)c(x) = -q(x)/p(x)[/tex]

Therefore, the self-adjoint second-order differential equation is transformed into the special Riccati equation:

(du/dx) + (a(x)u) + (b(x)u^2) + c(x) = 0

Now, let's apply this result to transform the self-adjoint equation:

(d/dx)(xy') + (1 - x)y = 0

We can rewrite this equation in terms of p(x) by setting p(x) = x:

(d/dx)(xy') + (1 - x)y = 0

Using the substitution v = p(x)y' = xy', we differentiate v with respect to x:

dv/dx = d/dx(xy')

Applying the product rule:

dv/dx = x(dy/dx) + y

Substituting v = xy' back into the original equation, we have:

(dv/dx) + (1 - x)y = x(dy/dx) + y + (1 - x)y = 0

Simplifying further:

x(dy/dx) + 2y - xy = 0

Comparing this equation with the general form of the Riccati equation:

[tex](du/dx) + a(x)u + b(x)u^2 + c(x) = 0[/tex]

We can identify the coefficients as:

a(x) = -x

b(x) = 0 (no u^2 term in the reduced equation)

c(x) = 2

Therefore, the self-adjoint equation is transformed into the Riccati equation:

(du/dx) - xu + 2 = 0

Applying this technique, the self-adjoint equation (d/dx)(xy') + (1 - x)y = 0 is transformed into the Riccati equation (du/dx) - xu + 2 = 0.

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For the first example, we are given a geometric sequence with the first three terms as 4, x, and x + 24.

To find the value of the car at a specific time, you need to calculate the depreciation for each year up to that time and subtract it from the initial value of £15,645.

In a geometric sequence, each term is found by multiplying the previous term by a constant called the common ratio.

Let's assume the common ratio is denoted by r.

Based on this information, we can write the following equations:

x = 4 × r,

x + 24 = x × r.

To find the possible values of x, we need to solve these equations simultaneously.

From the first equation, we can express r in terms of x: r = x/4.

Substituting this value of r into the second equation, we get:

x + 24 = (x/4) × x.

Simplifying this equation, we have:

4x + 96 = x².

Rearranging the equation, we get:

x² - 4x - 96 = 0.

Now we can solve this quadratic equation for x. Factoring or using the quadratic formula will yield the possible values of x.

For the second example, we are given that a car was purchased for £15,645 on 1st January 2021, and its value depreciates each year.

To determine the value of the car at a given time, we need to know the rate of depreciation.

Let's assume the rate of depreciation is d (expressed as a decimal).

The value of the car at the end of each year can be calculated as follows:

Year 1: £15,645 - d × £15,645,

Year 2: (£15,645 - d × £15,645) - d × (£15,645 - d × £15,645),

Year 3: [£15,645 - d × (£15,645 - d × £15,645)] - d × [£15,645 - d × (£15,645 - d × £15,645)],

and so on.

To find the value of the car at a specific time, you need to calculate the depreciation for each year up to that time and subtract it from the initial value of £15,645.

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The solution set of the system is {(x, y, z) = (2, 4, 2)}.

To solve the given system of equations using Cramer's Rule, we need to find the values of x, y, and z that satisfy all three equations simultaneously. Cramer's Rule involves calculating determinants to obtain the solution.

Find the determinant of the coefficient matrix (D):

Find the determinant of the x-column matrix (Dx):

Find the determinant of the y-column matrix (Dy):

Find the determinant of the z-column matrix (Dz):

Now, we can find the values of x, y, and z using the formulas:

Since the determinant of the coefficient matrix (D) is zero, Cramer's Rule cannot be applied to this system of equations. The system either has no solutions or infinitely many solutions. Therefore, the solution set of the system is empty, and there are no values of x, y, and z that satisfy all three equations simultaneously.

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The linear recurring sequence with characteristic polynomial 4 f(x) = x² + 5x³ + 2x² + 4 in F7[x] corresponds to a linear feedback shift register (LFSR). Its linear recurrence relation can be determined from the characteristic polynomial. The sequence is ultimately periodic, meaning it repeats after a certain number of terms. This is because the characteristic polynomial has a finite number of distinct roots in the field F7.

a) The corresponding LFSR (Linear Feedback Shift Register) for the given linear recurring sequence can be constructed by representing the characteristic polynomial as a feedback polynomial. The characteristic polynomial 4f(x) = x² + 5x³ + 2x² + 4 € F7[x] can be written as f(x) = x³ + 2x² + 4x + 4 € F7[x].

To draw the LFSR, we start with the shift register containing the initial values (S1, S2, S3) and the corresponding feedback connections represented by the coefficients of the polynomial. In this case, the LFSR would have three stages and the feedback connections would be as follows:

- The output of stage 1 is fed back to the input of stage 3.

- The output of stage 2 is fed back to the input of stage 1.

- The output of stage 3 is fed back to the input of stage 2.

b) In an ultimately periodic sequence, there exists a period and a pre-period. The period is the length of the repeating portion of the sequence, while the pre-period is the length of the non-repeating portion that leads to the repeating part.

The given linear recurring sequence is periodic because it satisfies the conditions for periodicity. The sequence is determined by a linear recurrence relation, which means each term is a function of the previous terms. As a result, the values of the sequence will eventually repeat after a certain number of terms. This repetition indicates the existence of a period.

Without computing the sequence explicitly, we can observe that the given sequence is ultimately periodic because it is generated by a linear recurrence relation with a finite number of terms. Once the sequence starts repeating, it will continue to repeat indefinitely. Therefore, the sequence is periodic.

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The matrix of the orthogonal projection on the two-dimensional subspace spanned by (1, 1, 1, 1) and (2, 0, -1, -1) in the standard basis of R4 is:

P =[tex]\left[\begin{array}{cccc}1/2&1/2&0&0\\1/2&1/2&0&0\\0&0&0&0\1&0&0&0&0\end{array}\right][/tex]

Here, we have,

To compute the matrix of an orthogonal projection on a two-dimensional subspace in R4, we need to find an orthonormal basis for that subspace first.

Here's the step-by-step process:

Step 1: Find the orthogonal complement of the given subspace.

Let's find a vector orthogonal to both (1, 1, 1, 1) and (2, 0, -1, -1).

Taking their cross product, we have:

(1, 1, 1, 1) × (2, 0, -1, -1) = (2, 2, -2, -2)

Step 2: Normalize the orthogonal vector.

Normalize the vector obtained in the previous step by dividing it by its length:

v = (2, 2, -2, -2) / √(16) = (1/2, 1/2, -1/2, -1/2)

Step 3: Find another orthogonal vector in the subspace.

Now, we need to find another vector in the subspace that is orthogonal to v.

We can choose any vector that is linearly independent of v.

Let's choose (1, 1, 1, 1).

Step 4: Normalize the second orthogonal vector.

Normalize the vector (1, 1, 1, 1) by dividing it by its length:

u = (1, 1, 1, 1) / 2 = (1/2, 1/2, 1/2, 1/2)

Step 5: Create an orthonormal basis for the subspace.

We now have two orthogonal vectors, v and u. To make them orthonormal, divide each vector by its length:

u' = u / ||u|| = (1/2, 1/2, 1/2, 1/2) / √(1/2) = (1/√2, 1/√2, 1/√2, 1/√2)

v' = v / ||v|| = (1/2, 1/2, -1/2, -1/2) /√(1/2) = (1/√2, 1/√2, -1/√2, -1/√2)

Step 6: Construct the projection matrix.

The projection matrix P can be constructed by taking the outer product of the orthonormal basis vectors:

P = u' * u'ⁿ + v' * v'ⁿ

Calculating this product, we have:

P = (1/√2, 1/√2, 1/√2, 1/√2) * (1/√2, 1/√2, 1/√2, 1/√2)ⁿ + (1/√2, 1/√2, -1/√2, -1/√2) * (1/√2, 1/√2, -1/√2, -1/√2)ⁿ

Simplifying this expression, we get:

P = (1/2, 1/2, 1/2, 1/2) * (1/2, 1/2, 1/2, 1/2) + (1/2, 1/2, -1/2, -1/2) * (1/2, 1/2, -1/2, -1/2)

P = (1/4, 1/4, 1/4, 1/4) + (1/4, 1/4, -1/4, -1/4)

P = (1/2, 1/2, 0, 0)

So, the matrix of the orthogonal projection on the two-dimensional subspace spanned by (1, 1, 1, 1) and (2, 0, -1, -1) in the standard basis of R4 is:

P =[tex]\left[\begin{array}{cccc}1/2&1/2&0&0\\1/2&1/2&0&0\\0&0&0&0\1&0&0&0&0\end{array}\right][/tex]

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To determine whether S = {(5, 4, 3), (0, 4, 3), (0, 0, 3)} is a basis for R3, we need to check if the vectors in S are linearly independent and if they span R3.

To check if the vectors in S are linearly independent, we can form a matrix with the vectors as its columns and perform row reduction. If the row-reduced echelon form of the matrix has a pivot in every row, then the vectors are linearly independent. If not, they are linearly dependent.

In this case, constructing the matrix and performing row reduction, we find that the row-reduced echelon form has a row of zeros. Therefore, the vectors in S are linearly dependent, and thus S is not a basis for R3.

Since S is not a basis for R3, we cannot write u = (15, 8, 15) as a linear combination of the vectors in S. The given expression, u = 3(5, 4, 3) - (0, 4, 3) + 3(0, 0, 3), does not yield the vector u = (15, 8, 15). Hence, the solution is IMPOSSIBLE.

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Therefore, the values of r in the general solution are r = 0 and r = 2.

To find the value of r in the general solution of the second-order differential equation 5y'' = 2y', we can rewrite the equation in standard form:

5y'' - 2y' = 0

Now, let's assume that the solution to this equation is of the form y(x) = c1eₓˣ + c2.

Taking the first and second derivatives of y(x), we have:

y'(x) = c1reˣ

y''(x) = c1r^2eˣ

Substituting these derivatives into the differential equation, we get:

5(c1r^2eˣ) - 2(c1reˣ) = 0

Simplifying the equation, we have:

c1(r² - 2r)eˣ = 0

For this equation to hold for all values of x, the coefficient of e^(rx) must be equal to zero:

r²- 2r = 0

Factoring out an r, we have:

r(r - 2) = 0

Setting each factor equal to zero, we get:

r = 0, r = 2

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"If it is not rainy, then I will go to the school" is the negation of "If it is rainy, then I will not go to the school".

To find the negation of a conditional statement, we need to reverse the direction of the implication and negate both the hypothesis and the conclusion.

The given statement is "If it is rainy, then I will not go to the school." Let's break it down:

Hypothesis: It is rainy

Conclusion: I will not go to the school

To negate this statement, we reverse the implication and negate both the hypothesis and the conclusion. The negation would be:

Negated Hypothesis: It is not rainy

Negated Conclusion: I will go to the school

So, the negation of "If it is rainy, then I will not go to the school" is "If it is not rainy, then I will go to the school." Therefore, the correct answer is option c) "If it is not rainy, then I will go to the school."

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There are three integer values of x (5, 6, and 7) that satisfy the condition √(2x) > 3 and √(2x) < 4.

To find the integer values of x that satisfy the condition √(2x) > 3 and √(2x) < 4, we need to consider the range of values for x that make the inequality true.

We start by isolating the square root expression:

3 < √(2x) < 4

To eliminate the square root, we can square both sides of the inequality:

3^2 < (√(2x))^2 < 4^2

9 < 2x < 16

Dividing the inequality by 2:

4.5 < x < 8

Now, we need to find the integer values of x that lie within this range. Since the condition asks for integer values, we can conclude that the possible values for x are 5, 6, and 7. Note that x cannot be equal to 4 or 8, as those values would make the inequality false.

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To solve this problem, let's go through each part step by step:

a) To compute the Pearson correlation coefficient, we need to calculate the covariance between the mother's education (X) and the father's education (Y), as well as the standard deviations of X and Y.

Given the data:

X (Mother's education): 10 8 10 7 15 4 9 6 N 12

Y (Father's education): 15 10 7 6 5 7 8 5 10 00

First, calculate the means of X and Y:

mean_X = (10 + 8 + 10 + 7 + 15 + 4 + 9 + 6 + N + 12) / 10 = (X + N) / 10

mean_Y = (15 + 10 + 7 + 6 + 5 + 7 + 8 + 5 + 10 + 0) / 10 = 6.8

Next, calculate the deviations from the mean for each data point:

deviations_X = X - mean_X

deviations_Y = Y - mean_Y

Compute the sum of the product of these deviations:

sum_of_product_deviations = Σ(deviations_X * deviations_Y)

Calculate the standard deviations of X and Y:

std_dev_X = √(Σ(deviations_X^2) / (n - 1))

std_dev_Y = √(Σ(deviations_Y^2) / (n - 1))

Finally, compute the Pearson correlation coefficient (r):

r = sum_of_product_deviations / (std_dev_X * std_dev_Y)

b) The coefficient of determination (r^2) is the square of the Pearson correlation coefficient. Therefore, r^2 = r^2.

c) To determine if there is a significant relationship between the mother's education and the father's education, we can conduct a two-tailed test using the t-distribution at a significance level of 0.05.

The null hypothesis (H0) is that there is no relationship between the mother's education and the father's education level.

The alternative hypothesis (H1) is that there is a significant relationship between the mother's education and the father's education level.

We can calculate the t-statistic using the formula:

t = r * √((n - 2) / (1 - r^2))

Next, we need to find the critical t-value for a two-tailed test with (n - 2) degrees of freedom and a significance level of 0.05. We can consult a t-table or use statistical software to find the critical value.

If the calculated t-statistic is greater than the critical t-value or less than the negative of the critical t-value, we reject the null hypothesis and conclude that there is a significant relationship between the mother's education and the father's education level.

d) To write the regression equation predicting the mother's educational level (X) from the father's education (Y), we can use the simple linear regression formula:

X = a + bY

where a is the intercept and b is the slope of the regression line.

To calculate the intercept and slope, we can use the following formulas:

b = r * (std_dev_X / std_dev_Y)

a = mean_X - b * mean_Y

e) To predict the mother's level of education (X) if the father has 15 years of education (Y = 15), we can substitute Y = 15 into the regression equation:

X = a + b * 15

Substitute the calculated values of a and b from part (d) into the equation and solve for x

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a) The directional derivative at (-1,2) in the direction of (2,1) is 0.

b) The surface described by function h is flat or constant in the direction of (2,1) at (-1,2).

c) There is no direction of fastest increase at (-1,2).

d) A contour plot for h can be generated using graphing software.

e) The contour plot visually represents the changing function values of h across different x and y values.

a) To find the directional derivative at (-1,2) in the direction of (2,1), we first compute the gradient of h(x, y), denoted as ∇h(x, y). The gradient is given by:

∇h(x, y) = (∂h/∂x, ∂h/∂y)

Taking partial derivatives, we have:

∂h/∂x = (2x) / (4 + x² + y²)

∂h/∂y = (2y) / (4 + x² + y²)

Evaluating these partial derivatives at (-1,2), we get:

∂h/∂x = (-2) / 5

∂h/∂y = (4) / 5

The directional derivative is then computed as the dot product of the gradient and the unit vector in the direction of (2,1). The unit vector is obtained by normalizing the direction vector:

u = (2,1) / √(2² + 1²) = (2,1) / √5 = (2/√5, 1/√5)

Finally, the directional derivative is:

D_u h(-1,2) = ∇h(-1,2) · u = (-2/5, 4/5) · (2/√5, 1/√5) = (-4/5√5) + (4/5√5) = 0

Therefore, the directional derivative at (-1,2) in the direction of (2,1) is 0.

b) The fact that the directional derivative is zero tells us that the surface described by the function h does not change in the direction of (2,1) at the point (-1,2). This means that the surface is flat or constant in that direction at that point.

c) To determine the direction of fastest increase at (-1,2), we look for the direction in which the directional derivative is maximized. Since the directional derivative is zero in this case, there is no direction of fastest increase at (-1,2).

e) A contour plot for h visually represents the level curves or contours of the function on a two-dimensional plane. The contour lines connect points with the same function value. By observing the contour plot, you can see how the function values change across different values of x and y. Areas with closely spaced contour lines indicate steep changes in the function value, while areas with widely spaced contour lines suggest slower changes. Additionally, contours that are close together suggest a steeper slope, while contours that are far apart indicate a flatter region of the surface.

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If f(x) is a linear function, it can be represented by the equation of a straight line in the form:

f(x) = mx + bwhere m is the slope of the line and b is the y-intercept.

Given that f(-5) = 3 and f(5) = 2, we can substitute these values into the equation to form a system of equations:

f(-5) = -5m + b = 3 ---- (1)

f(5) = 5m + b = 2 ---- (2)

To find the equation for f(x), we need to solve this system of equations for the values of m and

b.We can subtract equation (1) from equation (2) to eliminate the b term:5m + b - (-5m + b) = 2 - 3

5m + b + 5m - b = -1

10m = -1

m = -1/10

Substituting the value of m back into either equation (1) or (2) to solve for b:-5(-1/10) + b = 3

1/2 + b = 3

b = 3 - 1/2

b = 5/2

Therefore, the equation for f(x) is:

f(x) = (-1/10)x + 5/2

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The value of the integral ∫10 |x - 5| dx is 10.

Interpreting the integral in terms of areas, we can consider |x - 5| as a piecewise function that represents the absolute value of the difference between x and 5. The absolute value function ensures that the output is always positive or zero.

Since we are integrating over the interval [0, 10], we can split this interval into two regions: [0, 5] and [5, 10].

In the first region, where x is less than or equal to 5, |x - 5| simplifies to 5 - x. Integrating this function over the interval [0, 5] gives us an area of 10.

In the second region, where x is greater than 5, |x - 5| simplifies to x - 5. Integrating this function over the interval [5, 10] also gives us an area of 10.

Therefore, the total area under the curve |x - 5| over the interval [0, 10] is the sum of the areas in both regions, which is 10 + 10 = 20.

However, since the absolute value function ensures that the output is always positive or zero, the integral represents the signed area, which means areas below the x-axis are counted as negative. In this case, the integral evaluates to 10, representing the total net area between the curve and the x-axis over the interval [0, 10].

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Using implicit differentiation, the derivative dy/dx of the equation x² + y² = xy + 7 is found to be dy/dx = (y - x) / (y - 2x). Evaluating this at x = -3 and y = -2, we get dy/dx = 5/4.

To find dy/dx, we differentiate both sides of the equation x² + y² = xy + 7 with respect to x using the rules of implicit differentiation.

Differentiating x² + y² with respect to x gives 2x + 2yy' (using the chain rule), and differentiating xy + 7 with respect to x gives y + xy'.

Rearranging the terms, we have:

2x + 2yy' = y + xy'

Bringing the y' terms to one side and factoring out y - x, we get:

2x - y = (y - x)y'

Dividing both sides by y - x, we have:

y' = (2x - y) / (y - x)

Substituting x = -3 and y = -2 into the derivative expression, we get:

dy/dx = (y - x) / (y - 2x) = (-2 - (-3)) / (-2 - 2(-3)) = 5/4

Therefore, dy/dx evaluated at x = -3 and y = -2 is dy/dx = 5/4.


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The orthogonal complement of [1] is the set of all functions in V that satisfy this equation. This is a subspace of V that is spanned by the two functions x - 3/2 and x² - 3x + 15/2. The computation of (f, g) and || ƒ|| for f(x) = x + 2 and g(x)=x² - 2x - 3 is as follows:

Step by step answer:

1. To compute (f, g), use the given inner product: (f, g) = f(t)g(t)dt. Substitute f(x) = x + 2 and

g(x)=x² - 2x - 3:(f, g)

[tex]= ∫0¹ (x+2)(x²-2x-3)dx[/tex]

[tex]= ∫0¹ x³ - 2x² - 7x - 6dx[/tex]

[tex]= [-1/4 x^4 + 2/3 x^3 - 7/2 x^2 - 6x] |0¹[/tex]

[tex]= (-1/4 (1)^4 + 2/3 (1)^3 - 7/2 (1)^2 - 6(1)) - (-1/4 (0)^4 + 2/3 (0)^3 - 7/2 (0)^2 - 6(0))[/tex]

[tex]= -1/4 + 2/3 - 7/2 - 6= -41/12[/tex]

Therefore, (f, g) = -41/12.2.

To find || ƒ||, use the definition of the norm induced by the inner product: ||f|| = √(f, f).

Substitute f(x) = x + 2:||f||

= √(f, f)

= √∫0¹ (x+2)²dx

= √∫0¹ x² + 4x + 4dx

= √[1/3 x³ + 2x² + 4x] |0¹

= √[(1/3 (1)^3 + 2(1)^2 + 4(1)) - (1/3 (0)^3 + 2(0)^2 + 4(0))]

= √(11/3)

= √(33)/3

Thus, || ƒ|| = √(33)/3.3.

To find the orthogonal complement of the subspace of scalar polynomials, we first need to determine what that subspace is. The subspace of scalar polynomials is the span of the constant polynomial 1 on [0, 1], which is denoted by [1]. We need to find all functions in V that are orthogonal to all functions in [1].Let f(x) be any function in V that is orthogonal to all functions in [1]. Then we must have (f, 1) = 0 for all constant functions 1. This means that:∫0¹ f(x) dx = 0.

We know that the space of polynomials of degree ≤2 on [0, 1] has a basis consisting of 1, x, and x². Thus, any function in V can be written as:f(x) = a + bx + cx²for some constants a, b, and c. Since f(x) is orthogonal to 1, we must have (f, 1) = a∫0¹ 1dx + b∫0¹ xdx + c∫0¹ x²dx

= 0.

Substituting the integrals, we obtain: a + b/2 + c/3 = 0.This means that any function f(x) in V that is orthogonal to [1] must satisfy this equation. Thus, the orthogonal complement of [1] is the set of all functions in V that satisfy this equation. This is a subspace of V that is spanned by the two functions x - 3/2 and x² - 3x + 15/2.Another way to think about this is that the orthogonal complement of [1] is the space of all polynomials of degree ≤2 that have zero constant term. This is because any such polynomial can be written as the sum of a scalar polynomial (which is in [1]) and a function in the orthogonal complement.

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The vector z = [2 + 4i, -2i] can be written as a linear combination of z₁ and z₂ as,(z,z₁)z₁ + (z,z₂)z₂= (5 + 3i) [1 + i/2, 1 - i/2] + (-3√2 + i√2) [i/√2, -1/√2].

(a) Here, {z₁, z₂} is an orthonormal set in C².

We have given,

z₁ = [1 + i/2, 1 - i/2],z₂ = [i/√2, -1/√2].

Now, we need to show that {z₁, z₂} is an orthonormal set in C².As we know that, the inner product of two complex vectors v and w of dimension n is defined by the following formula:

(v,w) = ∑i=1nviwi^* where vi and wi are the i-th components of v and w, respectively, and wi^* is the complex conjugate of the i-th component of w.

(i) Inner product of z₁ and z₂ is

(1 + i/2).(i/√2) + (1 - i/2).(-1/√2)= i/(2√2) - i/(2√2) = 0

(ii) Magnitude of z₁ is∣z₁∣ = √((1 + i/2)² + (1 - i/2)²)= √(1 + 1/4 + i/2 + i/2 + 1 + 1/4)= √(3 + i)√((3 - i)/(3 - i))= √(10)/2

(iii) Magnitude of z₂ is∣z₂∣ = √((i/√2)² + (-1/√2)²)= √(1/2 + 1/2)= 1

(iv) Inner product of z₁ and z₁ is(1 + i/2).(1 - i/2) + (1 - i/2).(1 + i/2)= 1/4 + 1/4 + 1/4 + 1/4= 1

Therefore, {z₁, z₂} is an orthonormal set in C².

(b) Here, we are given z = [2 + 4i, -2i]and we need to write it as a linear combination of z₁ and z₂.

As we know that, we can write any vector z as a linear combination of orthonormal vectors z₁ and z₂ as,

z = (z,z₁)z₁ + (z,z₂)z₂where (z,z₁) = Inner product of z and z₁, and (z,z₂) = Inner product of z and z₂.

Now, let's calculate these inner products:

(z,z₁) = (z,[1 + i/2, 1 - i/2])

= (2 + 4i)(1 + i/2) + (-2i)(1 - i/2)

= 1/2 + 2i + 4i + 2 + i - 2i

= 5 + 3i(z,z₂)

= (z,[i/√2, -1/√2])

= (2 + 4i)(i/√2) + (-2i)(-1/√2)

= (2i - 4)(1/√2) + (2i/√2)

= -3√2 + i√2

Now, putting these values in the equation, we have z = (5 + 3i) [1 + i/2, 1 - i/2] + (-3√2 + i√2) [i/√2, -1/√2]

Thus, the vector z = [2 + 4i, -2i] can be written as a linear combination of z₁ and z₂ as,

(z,z₁)z₁ + (z,z₂)z₂

= (5 + 3i) [1 + i/2, 1 - i/2] + (-3√2 + i√2) [i/√2, -1/√2]

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The solution to the system of equations is: x = 11y = -48.

There are different methods to solve systems of linear equations but we will use the elimination method which involves the following steps: STEP 1: Multiply one or both of the equations by a suitable number so that one of the variables has the same coefficient in both equations. We have two equations:

24x + 3y = 792, 24x + (-b)y = 1464Multiplying the first equation by -1 will give us -24x - 3y = -792 and our equations now becomes:

-24x - 3y = -792 24x + (-b)y = 1464STEP 2: Add the two equations together. This eliminates one of the variables. We add the two equations together and simplify:

(-24x - 3y) + (24x - by) = (-792) + 1464Simplifying the left hand side, we have: -3y - by = 672Factorising y,

we have: y(-3 - b) = 672 y = -672/(3 + b)STEP 3: Substitute the value of y obtained into any one of the original equations and solve for the other variable.

Using the first equation:24x + 3y = 792 substituting y, we have:

24x + 3(-672/(3 + b)) = 792

Simplifying and solving for x, we have:24x - 224b/(3 + b) = 792

Multiplying both sides by (3 + b), we have:24x(3 + b) - 224b = 792(3 + b)72x + 24bx - 224b = 2376 + 792b

Collecting like terms: 72x + (24b - 224)b = 2376 + 792b72x + (24b² - 224b - 792)b = 2376Simplifying, we have:24b² - 224b - 792 = 0Dividing through by 8, we have:3b² - 28b - 99 = 0

Factoring the quadratic equation, we have:(3b + 9)(b - 11) = 0Therefore, b = -3 or b = 11Substituting b = -3, we have:y = -672/(3 - 3) = undefined which is not valid, hence b = 11

Therefore, y = -672/(3 + 11) = -48Therefore:x = (792 - 3y)/24 = (792 - 3(-48))/24 = 11 The solution to the system of equations is: x = 11y = -48.

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The point in spherical coordinates is now presented: (r, α, γ) = (4.216, - 18.434°, 46.506°)

In this problem we find the definition of a point in cylindrical coordinates, whose equivalent form is spherical coordinates must be found. We present the following definition:

(ρ · cos θ, ρ · sin θ, z) → (r, α, γ)

Where:

r = √(ρ² + z²)

γ = tan⁻¹ (ρ / z)

α = θ

Now we proceed to determine the spherical coordinates of the point: (ρ · cos θ = - 4, ρ · sin θ = 4 / 3, z = 4)

ρ = √[(- 4)² + (4 / 3)²]

ρ = 4.216

γ = tan⁻¹ (4.216 / 4)

γ = 46.506°

α = tan⁻¹ [- (4 / 3) / 4]

α = tan⁻¹ (- 1 / 3)

α = - 18.434°

(r, α, γ) = (4.216, - 18.434°, 46.506°)

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(a) Removing the absolute value, we get: i = ± e^(-5t + C1)

(b) the particular solution is: i_p = (22/3)sin(t)

(c) the particular solution for the given initial condition is:

i = (22/3)sin(t)

To solve the given differential equation, we'll first find the homogeneous solution and then the particular solution.

(a) Homogeneous Solution:

The homogeneous equation is given by:

L(di/dt) + RI = 0

Substituting the values L = 3 and R = 15, we have:

3(di/dt) + 15i = 0

Dividing by 3, we get:

(di/dt) + 5i = 0

This is a first-order linear homogeneous differential equation. We can solve it by separating variables and integrating:

(1/i) di = -5 dt

Integrating both sides, we get:

ln|i| = -5t + C1

Taking the exponential of both sides, we have:

|i| = e^(-5t + C1)

Removing the absolute value, we get:

i = ± e^(-5t + C1)

Now, let's find the particular solution.

(b) Particular Solution:

The particular solution is determined by the non-homogeneous term, which is E(t) = 110sin(t).

To find the particular solution, we assume i = A sin(t) and substitute it into the differential equation:

L(di/dt) + RI = E(t)

3(Acos(t)) + 15(Asin(t)) = 110sin(t)

Comparing coefficients, we get:

3Acos(t) + 15Asin(t) = 110sin(t)

Matching the terms on both sides, we have:

3A = 0 (to eliminate the cos(t) term)

15A = 110

Solving for A, we get:

A = 110/15 = 22/3

Therefore, the particular solution is:

i_p = (22/3)sin(t)

(c) Complete Solution:

The complete solution is the sum of the homogeneous and particular solutions:

i = i_h + i_p

i = ± e^(-5t + C1) + (22/3)sin(t)

Now, we can use the initial condition at t = 0, where I = 0, to determine the constant C1:

0 = ± e^(-5(0) + C1) + (22/3)sin(0)

0 = ± e^(C1) + 0

e^(C1) = 0

Since e^(C1) cannot be zero, we have:

± e^(C1) = 0

Therefore, the particular solution for the given initial condition is:

i = (22/3)sin(t)

(b) Finding the current after 15 minutes:

We need to find the value of i(t) after 15 minutes, which is t = 15 minutes = 15(60) seconds = 900 seconds.

Substituting t = 900 into the particular solution, we get:

i(900) = (22/3)sin(900)

Calculating sin(900), we find that sin(900) = 0.

Therefore, the current after 15 minutes is:

i(900) = (22/3)(0) = 0 Amps.

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